Stable Diffusion online AI

A prompt is a text description that guides the AI. It can include the subject, artistic style, lighting, mood, and medium — for example: 'a medieval knight in a dark forest, oil painting, dramatic lighting, cinematic'. The more specific your prompt, the closer the result to your vision. Keep it under 350 characters for best results.

A negative prompt tells the AI what to leave out of the image.

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A negative prompt tells the AI what NOT to show. Common examples: 'blurry, low quality, watermark, text, deformed hands'. Use these to remove typical AI mistakes.

Negative prompts boost image quality fast. Add weights — like '(blurry:1.4)' — to make them stronger.

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The seed is a number that initializes the random process behind image generation. Every image starts from noise, and the seed determines what that starting noise looks like. Leave it at -1 to get a random seed — and a different image — every time.

Fix the seed to reproduce the exact same image. This is useful for experimenting: keep the seed constant while changing only the prompt or settings to compare variations precisely. It is also how you recreate a favourite result later.

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LoRA stands for Low-Rank Adaptation. It is a small add-on file (typically 50–150 MB) that fine-tunes a base model for a specific style, character, or subject — without retraining the entire model.

Load a LoRA to generate images in a specific artistic style, reproduce a particular character, or match a reference concept. You can combine several LoRAs in one prompt using weights: <lora:my-style:0.8>. Thousands of free LoRAs are available on Civitai and Hugging Face.

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Running Stable Diffusion locally requires a recent NVIDIA GPU with at least 6 GB of VRAM. AMD GPUs work with ROCm drivers on Linux. Apple Silicon (M1/M2/M3) is also supported via the MPS backend. If you have no compatible GPU, our free online tool runs directly in your browser — no installation needed.

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You need at least 4 GB of VRAM to run Stable Diffusion. 6 GB is the recommended minimum for SDXL models. With 8 GB or more, you can generate larger images and use more features.

If you have no GPU, you can still run Stable Diffusion on CPU — but it is much slower. You can also use our free online demo with no GPU required.

Img2img lets you use an existing image as a starting point. You provide a source image and a prompt. Stable Diffusion then transforms the image based on your instructions.

Use it to change the style of a photo, fix a generated image, or refine a rough sketch. The denoising strength controls how much the AI changes the original — low keeps it close, high transforms it fully.

It depends on the model. The base Stable Diffusion models use the CreativeML Open RAIL-M license. This allows commercial use, but with some restrictions — you cannot use images to harm or deceive people.

Always check the license of the specific checkpoint you use. Some community models restrict commercial use. If you sell or publish AI-generated images, keep a record of the model and its license.

ControlNet is an extension for Stable Diffusion that gives you precise control over image composition. You provide a reference image — a pose, an edge map, or a depth map — and the AI follows that structure.

Use ControlNet to keep a person's exact pose, copy a scene layout, or trace a sketch precisely. It works alongside your prompt, so you control both the content and the structure of the image.

A checkpoint is a saved AI model file. It contains everything the model learned during training. You load a checkpoint to set the style and subject of your images.

Different checkpoints give very different results. Some are trained for realistic photos, others for anime or illustration. You can download free checkpoints from Hugging Face or Civitai.

ComfyUI is a free, node-based interface for Stable Diffusion. You build image workflows by connecting blocks called nodes. Each node does one step — loading a model, adding a prompt, or saving the image.

ComfyUI gives you full control over each step. It is faster and lighter than Automatic1111. Power users prefer it for complex workflows. You can share workflows as JSON files.

Flux is a text-to-image AI model by Black Forest Labs, founded by former Stability AI researchers. The latest version is FLUX.2, released in November 2025. It produces images up to 4MP with photorealistic output and multi-reference control.

FLUX.2 has four variants: Max (top quality), Pro (quality at full speed), Flex (precision and sharp text rendering), and Dev (open weights, free to deploy). FLUX.2 Dev can run locally in ComfyUI or Automatic1111 — the same tools as Stable Diffusion.

The latest release is Stable Diffusion 3.5 , released in October 2024. It produces sharper images with significantly better text rendering, richer colors, and improved fine details. It comes in three sizes: Large (8B parameters), Large Turbo for fast generation, and Medium — all runnable locally with moderate VRAM.

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