Stable Diffusion online AI
Artists Prompts Demo Chat GPT


Stable Diffusion

Artificial Intelligence Accessible to Everyone for free

AI robot painter
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What is Stable Diffusion ?

Stable Diffusion is a text-to-image Artificial Intelligence. It generates images from a simple description in natural language.

Developed by StabilityAI , it is open source. Unlike models like DALL-E or Midjourney, it can run on regular GPUs. SD community is very active. A lot of fine-tuned models and add-ons are available. You can even train a model with your own data !

Stable Diffusion uses part of the LAION Aesthetics dataset. It’s trained to make images that match text prompts. This model can make unique digital art. It can also do other things like animations and image manipulation.

You can see examples in our prompt gallery, Art Gallery and image browser.

Frequently Asked Questions


What is a prompt ?

A prompt is a sentence, usually 350 characters or less (75 tokens), that describes the image you want to generate.

Negative prompt is a simple way to specify what should not appear in generated image.


What is a negative prompt ?

A negative prompt is a list of words or concepts you want Stable Diffusion to avoid in the generated image. While the regular prompt describes what you want, the negative prompt acts as a filter — steering the model away from unwanted elements such as artifacts, deformations, or specific styles. Common entries include terms like blurry, bad anatomy, extra fingers, watermark, low quality, ugly.

Negative prompts work through classifier-free guidance: the model simultaneously moves toward your prompt and away from the negative prompt, with the CFG scale controlling how strongly both signals are applied. Note that Flux models largely make negative prompts obsolete — their improved architecture handles quality and anatomy natively, so prompt adherence alone is usually sufficient.


How does mentioning an artist influence the result ?

Adding an artist's name to your prompt biases the model toward their visual signature — brushwork, color palette, composition, lighting style, and level of detail. This works because Stable Diffusion was trained on large labeled datasets that include artworks attributed to specific artists. For example, adding Greg Rutkowski tends to produce highly detailed fantasy illustrations, while Monet shifts the output toward soft impressionist textures.

Not all artists carry the same weight in the model — some are strongly represented in the training data, others barely influence the output. The effect also varies between model versions: an artist well-recognized by SD 1.5 may be less prominent in SDXL or Flux. Explore hundreds of artist style comparisons on our dedicated artist reference page to find the styles that best match your creative vision.


What does seed mean ?

The seed is a number that starts the creation process. You don’t have to make this number yourself, it’s made randomly if you don’t choose one (usually by choosing -1).

But if you control the seed, you can make the same images again, try different settings, or change the prompt.


What is LoRA ?

LoRA means Low Rank Adaptation. It’s a set of small extensions that tweak base models.

You can use it to adjust Stable Diffusion for a certain style or topic. You can mix many LoRAs in one prompt with different weights. This opens up endless possibilities for creation.


Can I install Stable Diffusion locally ?

Yes ! a GPU with at least 6Gb memory (NVRAM) is required.

You may use customized models trained by the community, fine-tune results with LoRA and much more.


How to install Stable Diffusion locally ?

First, get the SDXL base model and refiner from Stability AI.
Next, make sure you have Pyhton 3.10 and Git installed.
Then, download and set up the webUI from Automatic1111.
Put the base and refiner models in this folder: models/Stable-diffusion under the webUI directory.
Find webui.bat in the main webUI folder and double-click it. Wait for it to install.
Finally, open your web browser and go to your local host on port 8760 to use Stable Diffusion.

What is ComfyUI ?

ComfyUI is a node-based graphical interface for Stable Diffusion. Rather than a traditional form, it lets you build image generation pipelines visually by connecting nodes — each node handles one specific step: loading a model, encoding a prompt, sampling, or decoding the image. This makes the entire process transparent and fully customizable.

Compared to Automatic1111, ComfyUI offers more granular control over the generation pipeline and is generally faster, at the cost of a steeper learning curve. It supports all major models: SD 1.5, SDXL, Stable Diffusion 3, and Flux. The ecosystem is extended via ComfyUI Manager, which lets you install hundreds of community custom nodes directly from the interface. Download it from the official ComfyUI GitHub repository.


What is Flux ?

Flux is a state-of-the-art text-to-image model released in 2024 by Black Forest Labs, founded by former Stability AI researchers. It significantly outperforms SDXL in prompt adherence, photorealism, and — notably — the ability to render legible text inside generated images, a long-standing weakness of previous diffusion models.

Flux comes in three variants: FLUX.1-schnell (fastest, Apache 2.0 license), FLUX.1-dev (higher quality, open weights for non-commercial use), and FLUX.1-pro (best quality, API only). The open variants are available on Hugging Face and run locally via ComfyUI or Automatic1111 with the appropriate extension.


What is the newest version of Stable Diffusion ?

On October 23, 2024, Stability AI released Stable Diffusion 3.5 . It is their most capable text-to-image model with great impovement with spelling abilities, performance and quality.


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© 2024 • All images are free to use • Developed by @anthony for Stable Diffusion online AI